StableDiffusion
So I work at a home improvement store, and one of my co-workers does some contracting work on the side. He is trying to encourage one of his neighbors to put some simple small-park kind of stuff on a plot of land he owns so that he (my co-worker) can pick up some extra business installing it. He's seen me messing around with Stable Diffusion on some web apps at work on my down time, and he asked if it was possible if I could take a photo of the site and use AI to insert some of these elements into it so that he could show it to this potential client and maybe sell it to him that way. "Sure," I said, thinking to myself, 'I can just use inpainting to blend this stuff into the image pretty seamlessly. Easy-peasy.' It took me almost a full day of on-again, off-again work to get a picnic table I could live with. But I CANNOT get any model, any prompt, anything to make a swing set that I can live with. I've been pecking away at this problem for several days now, and every single attempt at a swing set has resulted in something that is mangled, twisted, or some terrible hybrid of OTHER playground equipment that my co-worker definitely doesn't want in the scene. At this point I'm just working on it for the challenge, but I admit that I'm stumped. Short of training my own Lora, does anyone have any advice on how to make a coherent swing set to bring into this image? \>\_\< Yes, this is a silly problem to have, I admit that, but I've also learned a great deal about how Stable Diffusion 'thinks' in the last few days, so I consider it a learning experience. Still, does anyone else have any ideas?
cross-posted from: https://lemmy.dbzer0.com/post/1352467 > These have all been created with the base SDXL release + refiner in Automatic1111's using the extension to add the "refiner" found (https://github.com/wcde/sd-webui-refiner) > > 1216x832 resolution with Hi-Res Fix at 2x using the Siax-200k upscaler and DPM 2 a Karras sampler. > > Here's the dynamic prompt. You need to install the dynamic prompt extension from (https://github.com/adieyal/sd-dynamic-prompts) > > This is the prompt I used > > ``` {natural | forest | mountain | beach | desert| rural | urban | city | sci-fi | space | dystopian| futuristic}, > > {cinematic shot | establishing shot | intimate scene | sweeping grandeur},(21:9 aspect ratio) > > {(ultrawide panoramic cinematic aspect ratio)} > > {Blender | Photoshop|Octane render|Unreal Engine 5|8K RAW Photo} , > > {2-4$$lone twisted tree | winding river| mountain peak| crumbling ruins| abandoned cabin|wooden fence | dramatic cliffs | stormy sea | rolling thunder | howling wind | foggy moor | charred forest | broken-down cart| towering dunes | parched canyon | bone-strewn pit | petrified woods| wrecked galleon | beast's den| majestic waterfall | calm lake | moonlit trail | moss-covered stones | misty vale |ravaged battlefield | derelict mill} > > {cirrus clouds |stormy sky|cumulus clouds|stratus clouds|nimbostratus clouds|cumulonimbus clouds} > > {clear | atmospheric fog | mist | haze | pollution| dust |smoke |atmospheric halo| sun dogs | moon dogs | sun pillars | circumzenithal arcs|circumhorizontal arcs}, > > {abstracted | concept art| Hyperrealistic| stylized| fantasy| impressionistic | photo| realistic } > > (16K, 32bit color, HDR, masterpiece, ultra quality) > > {brutalist | minimalist| whimsical| retro futurist} > > {muted tones | vibrant hues} > > {warm sunset tones |cool muted blues | colors} > > {natural | warm| dramatic } > > {god rays | sun beams | crepuscular rays| antisolar rays | volumetric light | light pillars | sun pillars | moon pillars}, > > {dawn | sunset| night} {clear | overcast | fog} > > {winter | spring | summer | autumn} > > { volumetric shadows | volumetric ambiance | aerial perspective | depth fog}, > > in the style of > {1-2$$Dylan Furst |Ash Thorp | Simon Stålenhag | Bob Ross| Ralph McQuarrie | Syd Mead| Moebius| Daarken| Felix Yoon| Gustave Doré| Arnold Böcklin| William Blake | Frank > Frazetta| John Constable |J.C. Dahl } > and > {1-2$$James Gurney | Craig Mullins| Android Jones |Justin Maller | John Berkey| Roger Dean| Rodney Matthews | Chris Foss| Nicolas Roeg | Geoffrey Hayes | John Harris| Dinotopia| Jon Foster| Brom| Brian Froud | Alan Lee}, > ``` > >  > >  > >  > >  > >  > >  > >  > >  > >  > > 
Tried my hand at cars. I used the Juggernaut model. The prompt uses regional prompter extension and looks like this: > > > at night city, summer, sweaty ADDCOMM > RAW photo, Nikon Z6 II Mirrorless Camera, hyper realism, extremely detailed, 8k uhd, dslr, soft lighting, high quality, film grain ADDBASE > ADDROW > ADDCOL > Aston Martin zeekars, hotrod, LED, [lora:zeekars:.7](lora:zeekars:.7) ADDROW > Negative prompt: JuggernautNegative-neg, transition of shapes, blurry, (((numbers on door))), duplicates, close range angle, crop, illustration, drawing, painting, sketching, render, artwork, 3d, cgi, logo, text, letterbox, 3D, render, video game, anime, cartoon, sketch, caption, subtitle, signature, watermark, username, artist name > > I use SD upscale and ControlNet Tile for upscaling with 4x\_NMKD-Siax\_200k. zeekars is a Lora. I hope SDXL will improve in all the small details for mechanical, geometrical forms. At first, the image looks great but when you look closely, panels separation, headlight, tyres look unrealistic...We will see tomorrow!
Just wanted to put this out there for anyone else who was in the same position, as I'd spent some time banging on this to find a functioning combination and would have appreciated having had success reports myself. Running Debian Trixie, current as of July 22, 2023. I see 512x512 speeds of about 2.2 it/s, which is significantly slower than an lower-end Nvidia card than I'd used and significantly (about 1/8th the speed) that other people [have reported](https://old.reddit.com/r/StableDiffusion/comments/139lz6c/does_it_mean_stable_diffusion_performance_on/) getting the same XT 7900 XTX card running at on Linux), so there is probably more work for me to do, but it's definitely running on the GPU and is much faster than running on the CPU, so I know that this combination (vanilla system Python, vanilla system drivers, torch nightly in a venv) does at least *work*, which was something that I'd been unsure of up until now. Running on the host, no Docker containers. Using a venv. Automatic1111 web UI, in-repository drivers, 2.1.0.dev20230715+rocm5.5 torch via pip installed in a venv, standard system Python 3.11 (i.e. did not need to set up Python 3.8, as I've seen some people do). Needs the `non-free-firmware` apt repo component enabled; I have `firmware-amd-graphics-20230515-3`. Rocm 5.6 is out as of this writing from AMD, but Debian Trixie presently only has 5.5 packaged and in the repos. I did need to install `libstdc++-13-dev` -- only `libstdc++-12-dev` being installed caused Automatic1111 to bail out with an error in not being able to find a `limits` C++ header when building some C++ code at runtime; some users had run into a similar error and resolved it by installing `libstdc++-12-dev`, which was a bit confusing. I have both clang and g++ installed. I am not terribly familiar with the AMD ROCM stack, but my understanding is that part of it (libamdhip64?) performs some compilation at runtime; it apparently remembers the binaries it has compiled, as if I removed libstdc++-13-dev after a successful run, it continued to work. The user running the Automatic1111 frontend needed to be added to the `render` and `video` groups to have access to the requisite device files. I did not need to have `HSA_OVERRIDE_GFX_VERSION` set. As for options being passed in `COMMAND_ARGS`, just `--medvram` and `--api`. `--xformers` does not work with AMD cards; Stable Diffusion (or Automatic1111, unsure about responsibility in the stack) apparently just ignores it there; passing it doesn't break anything. Some `--opt-sdp` options, like `--opt-sdp-attention`, cause dramatic slowdown, I assume causing the generation to run on the CPU instead of the GPU. I'd suggest to anyone trying to get a similar environment running to not start including optimization flags until they have things working without them; this had complicated things for me. I see 2.59 it/s, so something like 20% higher performance, without `--medvram` being passed to `COMMANDLINE_ARGS`. I have not done extensive testing to see whether any issues show up elsewhere with Stable Diffusion.
Hi, For a media project, I need to create dark fantasy themed backgrounds. It will be heavily inpainted to meet each background needs. I'm looking for models, loras, styles, examples, tutorials, etc. Anything you can think around dark fantasy is valuable, please feel free to suggest linked or related subjects or resources. Thanks for help in advance.
I am really enjoying the model [aZovyaPhotoreal\_v2](https://civitai.com/models/57319/a-zovya-photoreal) for creating sci-fi space like images. Generating data \-- parameters long shot scenic professional photograph of full shot, people, one woman, auburn hair, blue eyes , pilot officers, on the command bridge of a giant space ship, overlooking a nebula, robotech, crew, detailed skin texture, (blush:0.5), (goosebumps:0.5), subsurface scattering Negative prompt: NSFW, Cleavage, Pubic Hair, Nudity, Naked, Au naturel, Watermark, Text, censored, deformed, bad anatomy, disfigured, poorly drawn face, mutated, extra limb, ugly, poorly drawn hands, missing limb, floating limbs, disconnected limbs, disconnected head, malformed hands, long neck, mutated hands and fingers, bad hands, missing fingers, cropped, worst quality, low quality, mutation, poorly drawn, huge calf, bad hands, fused hand, missing hand, disappearing arms, disappearing thigh, disappearing calf, disappearing legs, missing fingers, fused fingers, abnormal eye proportion, Abnormal hands, abnormal legs, abnormal feet, abnormal fingers Steps: 25, Sampler: DPM2 a Karras, CFG scale: 7, Seed: 3402089243, Face restoration: CodeFormer, Size: 768x512, Model hash: 5594efef1c, Model: aZovyaPhotoreal\_v2, Denoising strength: 0.4, Hires upscale: 1.75, Hires upscaler: None, Version: v1.4.0 postprocessing Postprocess upscale by: 3, Postprocess upscaler: Nearest, Postprocess upscaler 2: R-ESRGAN 4x+, CodeFormer visibility: 0.324, CodeFormer weight: 0.296 extras Postprocess upscale by: 3, Postprocess upscaler: Nearest, Postprocess upscaler 2: R-ESRGAN 4x+, CodeFormer visibility: 0.324, CodeFormer weight: 0.296
jordanwages.com
Full Post Text --- **In the spirit of full disclosure**, the content of this post is heavily cribbed from [this post on Reddit](https://www.reddit.com/r/StableDiffusion/comments/zyi24j/how_to_turn_any_model_into_an_inpainting_model/). However, [as we've seen](https://www.msn.com/en-us/news/technology/some-redditors-say-theyre-walking-away-after-apollo-app-shuts-down/ar-AA1dillN), the Internet is not forever. It is entirely possible that a wealth of knowledge could be lost at any time due to any number of reasons. Because I have found this particular post so helpful and find myself coming back to it over and over, I thought it would be appropriate to share this method of creating an inpainting model from any custom stable diffusion model. Inpainting models, like the name suggests, are specialized models that excel at "filling in" or replacing sections of an image. They're especially good at decoding what a section of an image should look like based on the section of image that already exists. This is useful when you're generating images and only a small section needs to be corrected, or if you're trying to add something specific to an image that exists. So how is this done? With a model merge. [Automatic1111](https://github.com/AUTOMATIC1111/stable-diffusion-webui) has an excellent model merging tool that we'll use. Let's assume that you have a custom model called `my-awesome-model.ckpt` that is based on `stable-diffusion-1.5`. In the A1111 Checkpoint Merger interface, follow these steps: 1. Set "Primary model (A)" to `stable-diffusion-1.5-inpainting.ckpt`. 2. Set "Secondary model (B)" to `my-awesome-model.ckpt`. 3. Set "Tertiary model (C)" to `stable-diffusion-1.5.ckpt`. 4. Set "Multiplier (M)" to `1`. 5. Set "Interpolation Method" to `Add difference`. 6. Give your model a name in the "Custom Name" field, such as `my-awesome-model-inpainting.ckpt`. * Adding "-inpainting" will signal to A1111 that the model is an inpainting model. This is useful for extensions such as [openOutpaint](https://github.com/zero01101/openOutpaint). Also, it's just a good idea to properly label your models. Because we know you're a degenerate that has hundreds of custom waifu models downloaded from [CivitAI](https://civitai.com/). 7. Click `Merge`. And bazinga! You have your own custom inpainting stable diffusion model. Thanks again to [/u/MindInTheDigits](https://www.reddit.com/user/MindInTheDigits) for sharing the process.
mastodon.jordanwages.com
This is a Mastodon thread I created featuring devices and computers that never were, but could have been. I think my favorite might be the [HD Laserdisc player called the MOID](https://mastodon.jordanwages.com/@wagesj45/110686626102818960). [https://mastodon.jordanwages.com/system/media\_attachments/files/110/686/620/475/187/344/original/95ff43d76ba41dd6.jpg](https://mastodon.jordanwages.com/system/media_attachments/files/110/686/620/475/187/344/original/95ff43d76ba41dd6.jpg)
Learning to inpaint and messing around with a [steampunk lycoris](https://civitai.com/models/87257)
The reference family of preprocessors in ControlNet allow you to quickly transfer a style from an image in txt2img or img2img mode. I will demonstrate how to use it in txt2img. * download this [image](https://pixabay.com/illustrations/digital-illustration-jesus-christ-8052065/). I chose the 853x1280 size. We will use it as our ControlNet Reference image. * setup your ControlNet tab like this. [ControlNet setup](https://i.imgur.com/wKcCX5K.jpg) If you wonder what is the difference between the reference preprocessors, you can find more info [here](https://github.com/Mikubill/sd-webui-controlnet/discussions/1280) \-add a prompt close to the style of your image. **Very important, set the CFG to a low value (3,4)**. Otherwise, your result will look very dark and overcooked. This is my settings: [text2img settings](https://i.imgur.com/1GtP7FH.jpg) Try to keep the aspect ratio of the original image but do not hesitate to experiment. * Results: [original image](https://i.imgur.com/OF7sXv6.png) [Albert](https://i.imgur.com/qi3gY5a.png) [Scarlett](https://i.imgur.com/Kowj0U5.png)
I thought I would share my experience this evening with the group here seeing that I still excited as hell for getting this hodgepodge to work at all. I have been playing with the machineMl version of stable diffusion on Windows for a while now (we won't go into the reasons why, but I have the 6800XT, which is not well suited to this use case) Automatic11111 on MachineML is dog slow, but I got some interesting results. So today I cleared out an old SSD, wired it up and installed a clean ubuntu. Following this [guide](https://www.youtube.com/watch?v=2XYbtfns1BU) I managed to get rocM running and the difference is like chalk and cheese. Or rather Impulse and Warp drive. Totally crazy! So for you AMD Radeon users out there. There is hope. It is in Linux, but it is there.
AI Art News & Culture by the people, for the people. Click to read Stable Digest, by Stable Foundation, a Substack publication with tens of thousands of readers.
Random illustration for a story idea; the story idea may not pan out, but I was proud of how the art came out! Steps: 40, Sampler: Euler a, CFG scale: 10, Seed: 900795974, Size: 1536x1536, Model hash: 25ba966c5d, Model: aZovyaRPGArtistTools\_v3, Denoising strength: 0.3, Clip skip: 2, Token merging ratio: 0.5, Ultimate SD upscale upscaler: 4x-AnimeSharp, Ultimate SD upscale tile\_width: 512, Ultimate SD upscale tile\_height: 512, Ultimate SD upscale mask\_blur: 8, Ultimate SD upscale padding: 32, Version: v1.3.2
I hope it’s the open release of SDXL - the beta on the Stable Diffusion discord is getting pretty impressive. In any case, I’d like a tea serving drone 😁.
Well, it's a mutated Arnold Schwarzenegger riding a dinosaur on the moon with a (dinosaur-)moon in the background! Not sure how effective the headlight is, but I guess it is how it is! I actually didn't keep the prompts, but it was not that difficult. It was made with the recent Dreamshaper model.
I [made](https://zavant.ai/creations/4112e8fe-59c4-457b-bef8-2a21e3183cf2) a large scale working QR Code scene that points to kbin.social. I used the new QR Code Control for SD 1.5 (released [here](https://huggingface.co/DionTimmer/controlnet_qrcode), and multiple rounds of upscale in img2img using it, controlnet tile, and the t2ia color control, starting from a QR code I made with [this QR code generator](https://quickqr.art/).
lemmy.dbzer0.com
Hey peeps, someone posted about this magazine to the stable_diffusion community in my instance. I just thought I'd inform you that the official /r/stablediffusion failover is over here: * [kbin link](https://kbin.social/m/stable_diffusion@lemmy.dbzer0.com) * [lemmy relative link](/c/stable_diffusion@lemmy.dbzer0.com)
cross-posted from: https://lemmy.ml/post/1153505 > Result of just playing around
StableDiffusion
!StableDiffusion@kbin.socialFor discussions around Stable Diffusion, a text-to-image generative AI model. Share your generated pictures, discuss the various UI and extensions, share news about releases, bring tutorials and more!